Stable diffusion image to image guide. This will let you run the model from your PC.

We will use AUTOMATIC1111 Stable Diffusion WebUI. A decoder, which turns the final 64x64 latent patch into a higher-resolution 512x512 image. Users can generate NSFW images by modifying Stable Diffusion models, using GPUs, or a Google Colab Pro subscription to bypass the default content filters. Here's a step-by-step guide: Load your images: Import your input images into the Img2Img model, ensuring they're properly preprocessed and compatible with the model architecture. Principle of Diffusion models (sampling, learning) Diffusion for Images – UNet architecture. Putting them all together, the prompt is. Feb 16, 2023 · Key Takeaways. 4. Create beautiful art using stable diffusion ONLINE for free. 5. As good as DALL-E (especially the new DALL-E 3) and MidJourney are, Stable Diffusion probably ranks among the best AI image generators. This process is repeated a dozen times. It uses text prompts as the conditioning to steer image generation so that you generate images that match the text prompt. Click the Send to Inpaint icon below the image to send the image to img2img > inpainting. White pixels are inpainted and black pixels are preserved. Images generated by Stable Diffusion based on the prompt we’ve provided. Choose the v1. This mask will indicate the regions where the Stable Diffusion model should regenerate the image. The same seed and the same prompt given to the same version of Stable Diffusion will output the same image every time. Jul 7, 2024 · You can use ControlNet along with any Stable Diffusion models. Mar 19, 2024 · Additional details – These are keywords that are more like sweeteners,e. The predicted noise is subtracted from the image. They are used to generate images, most commonly as text to image models: you give it a text prompt, and it returns an image. It is created by Stability AI. Aug 9, 2023 · Stable Diffusion checkpoint est le modèle que vous utilisez, ici Stable Diffusion 1. This approach is chosen for its efficiency, as operating in the smaller latent space results in a considerably faster process. The resulting image keeps the colors and layout of the original picture, letting us add a personal touch with text to turn simple sketches into awesome artwork. The noise predictor then estimates the noise of the image. The Web UI offers various features, including generating images from text prompts (txt2img), image-to-image processing (img2img This action will initialize the model and provide you with a link to the web interface where you can interact with Stable Diffusion to generate images. The extra Mar 19, 2024 · In AUTOMATIC1111 GUI, Select the img2img tab and select the Inpaint sub-tab. Check out the AUTOMATIC1111 Guide if you are new to AUTOMATIC1111. A diffusion model, which repeatedly "denoises" a 64x64 latent image patch. The super resolution component of the model (which upsamples the output images from 64 x 64 up to 1024 x 1024) is also fine-tuned, using the subject’s images exclusively. With its 860M UNet and 123M text encoder, the How to use Stable Diffusion on a Mac. Send to During training, instead of generating noisy images, the model generates tensor in latent space. This is an image generation application based on the Stable Diffusion model, capable of producing high-quality and diverse image content. The model was pretrained on 256x256 images and then finetuned on 512x512 images. In other words, the following relationship is fixed: seed + prompt = image Stable Diffusion is a Latent Diffusion model developed by researchers from the Machine Vision and Learning group at LMU Munich, a. Strength. Create a dedicated folder, you can call it stable diffusion (or any other name you prefer). You should set it to ‘ Whole Picture ’ as the inpaint result matches better with the overall image. Every image you generate in this web interface will be kept in the following directory (unless you modify the settings for someplace different): C:\stable-diffusion-webui\outputs\txt2img-images. Stable diffusion pipelines Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. A widgets-based interactive notebook for Google Colab that lets users generate AI images from prompts (Text2Image) using Stable Diffusion (by Stability AI, Runway & CompVis). This model is trained in the latent space of the autoencoder. Unlike the other two, it is completely free to use. adding some interesting details. Stable Diffusion NSFW refers to using the Stable Diffusion AI art generator to create not safe for work images that contain nudity, adult content, or explicit material. The most popular image-to-image models are Stable Diffusion v1. It determines how much of your original image will be changed to match the given prompt. This is the tile size to be used for SD upscale. The green recycle button will populate the field with the seed number used in In this tutorial I'll go through everything to get you started with #stablediffusion from installation to finished image. Set both the image width and height to 512. Use the paintbrush tool to create a mask. Or, if you've just generated an image you want to upscale, click "Send to Extras" and you'll be taken to there with the image in place for upscaling. A Quick Start Guide to help you decide how to use Stable Diffusion. It’s a free and popular choice. For anime images, it is common to adjust Clip Skip and VAE settings based on the model you use. This model uses a frozen CLIP ViT-L/14 text May 12, 2023 · Stable diffusion is an open-source machine learning model built to generate images from text, modify the images based on text, or fill in details on low-resolution or low-detail images. It is convenient to enable them in Quick Settings. Nov 2, 2022 · The image generator goes through two stages: 1- Image information creator. With your images prepared and settings configured, it's time to run the stable diffusion process using Img2Img. Fix details with inpainting. Let’s make some images Apr 3, 2024 · Here in our prompt, I used “3D Rendering” as my medium. A diffusion model is a type of generative model that's trained to produce stuff. Max: 512x512 px (auto-resized) Jan 27, 2024 · Adding Conditional Control to Text-to-Image Diffusion Models. “A serene landscape photograph of a tranquil lake reflecting the rugged peaks of the Rockies, surrounded by dense pine forests. 0 model, SSD-1B boasts significant improvements: it's 50% smaller in size and 60% 🤗 Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. For example, if you type in a cute Sep 6, 2023 · Stable Diffusion AI is a type of AI model that uses a process called diffusion to generate images. 5 model. Aug 19, 2023 · Mastering the art of prompting in Stable Diffusion is a journey that requires practice, experimentation, and a keen understanding of the nuances that shape the image generation process. Aug 22, 2022 · Go back to the create → Stable page again if you’re not still there, and right at the top of the page, activate the “Show advanced options” switch. Apr 3, 2024 · Step 3: Fine-tuning and Personal Touches. Iterate if necessary: If the results are not satisfactory, adjust the filter parameters or try a different filter. Another flabbergasting feature of Stable Diffusion is it’s ability to paint in new features to a pre-existing image…or at least to try! Dec 26, 2023 · Step 1: Upload the image to AUTOMATIC1111. Software setup. But they can also be used for inpainting and outpainting, image-to-image (img2img) and a lot more. It's trained on 512x512 images from a subset of the LAION-5B database. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of Feb 17, 2024 · What is Stable Diffusion WebUI (AUTOMATIC1111) Why AUTOMATIC1111 Is Popular Installing Stable Diffusion WebUI on Windows and Mac Installing AUTOMATIC1111 on Windows Installing AUTOMATIC1111 on Apple Mac Getting Started with the txt2img Tab Setting Up Your Model Crafting the Perfect Prompt Negative Prompts Fiddling with Image Size Batch Settings Guiding Your Model with CFG Scale Seed and May 17, 2023 · Image Generation Implementation: Choose the AI implementation that's used for image generation. Otherwise, you can drag-and-drop your image into the Extras Mar 14, 2023 · The default setting for Seed is -1, which means that Stable Diffusion will pull a random seed number to generate images off of your prompt. Stable Diffusion v3 introduces a significant upgrade from v2 by shifting from a U-Net architecture to an advanced diffusion transformer architecture. Mar 29, 2024 · Segmind Stable Diffusion-1B, a diffusion-based text-to-image model, is part of a Segmind's distillation series, setting a new benchmark in image generation speed, especially for high-resolution images of 1024x1024 pixels. Mar 28, 2023 · The sampler is responsible for carrying out the denoising steps. cinematic lighting, rim lighting. Lighting – Controling light is important for a good image. Open the provided link in a new tab to access the Stable Diffusion web interface. So once you find a relevant image, you can click on it to see the prompt. This approach aims to align with our core values and democratize access, providing users with a variety of options for scalability and quality to best meet their creative needs. Step 2: Select an inpainting model. It's one of the most widely used text-to-image AI models, and it offers many great benefits. Stable Diffusion XL is the latest iteration of the popular text-to-image generation model, offering impressive results. Only Masked Padding: The padding area of the mask. Upload the image to the inpainting canvas. You Mar 19, 2024 · Stable Diffusion Models: a beginner’s guide. This model is designed to convert textual descriptions into high-resolution, detailed images, which allows you to generate better NSFW or Porn pictures. g. Sep 21, 2022 · Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. ckpt command from the v1. Écrivez le prompt que vous voulez et une image apparaîtra. This will let you run the model from your PC. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. What is Img2Img in Stable Diffusion. Be detailed and specific in your prompt. Understanding prompts – Word as vectors, CLIP. 5 and 0. Feb 22, 2024 · The Stable Diffusion 3 suite of models currently ranges from 800M to 8B parameters. emaonly. Image2Image Pipeline for Stable Diffusion using 🧨 Diffusers. May 1, 2024 · Sampling steps are the number of iterations Stable Diffusion must go through, in order to generate a recognizable image out of random noise. You need to make sure there is at least 10 GB of free space. 2. Description. A 1-click Colab notebook I used for generating the site's images. Considering its roots, it is expected to become the new Mar 4, 2024 · Step Two: Find some Checkpoints. 3D rendering. The results from the Stable Diffusion and Kandinsky models vary due to their architecture differences and training process; you can generally expect SDXL to produce higher quality images than Stable Diffusion v1. After applying stable diffusion techniques with img2img, it's important to Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. When inpainting, setting the prompt strength to 1 will . It’s time to add your personal touches and make the image truly yours. It is suitable for various creative tasks, where you can simply choose or input the appropriate prompt to instantly generate images. Nov 24, 2023 · Image-to-image starts with an image you specify and then adds noise. Feb 18, 2024 · Inpaint Area: This lets you decide whether you want the inpainting to use the entire image as a reference or just the masked area. You should now be on the img2img page and Inpaint tab. Mar 16, 2024 · For local AI image generation, it’s hard to beat Stable Diffusion. If you are an AI art enthusiast, you might hear about Midjourney, DALL-E, or anything Jan 19, 2024 · Img2img in stable diffusion, also known as image-to-image, is a method that creates new AI images from a picture and a text prompt. Convert to landscape size. Stable Diffusion Models, or checkpoint models, are pre-trained Stable Diffusion weights for generating a particular style of images. This component is the secret sauce of Stable Diffusion. Diffusion in latent space – AutoEncoderKL. Basic implementations of Stable Diffusion can accept three inputs: Prompts. It begins by training a model on a dataset of images. Let words modulate diffusion – Conditional Diffusion, Cross Attention. Stable Diffusion is a collection of open-source models by Stability AI. Stable Diffusion is named that way because it's a latent diffusion model. All of Stable Diffusion's upscaling tools are located in the "Extras" tab, so click it to open the upscaling menu. This notebook shows how to create a custom diffusers pipeline for text-guided image-to-image generation with Stable Diffusion model using 🤗 Hugging Face 🧨 Diffusers library. They guide Stable Diffusion by defining the regions to be filled or preserved. The textual inversion tab within the web UI serves as In image editing, inpainting is a process of restoring missing parts of pictures. Running Stable Diffusion Locally Jul 31, 2023 · Check out the Quick Start Guide if you are new to Stable Diffusion. Stable Diffusion 3is an open-source diffusion model, the long waited upgrade to SDXL. May 16, 2024 · From there, select the 'inpaint' option and upload your image to initiate the process. Init image. Also, SDXL is an evolution of the previous Stable Diffusion models, offering significant We would like to show you a description here but the site won’t allow us. Stable Diffusion cannot read your mind. It can create images in variety of aspect ratios without any problems. The default settings are pretty good. color – The color scheme of the image. Here’s where your vision meets technology: enter a prompt that encapsulates the Jan 30, 2024 · Stable Diffusion is a text-to-image model, powered by AI, that uses deep learning to generate high-quality images from text. Oct 9, 2023 · The subject’s images are fitted alongside images from the subject’s class, which are first generated using the same Stable Diffusion model. Prompt string along with the model and seed number. Powered By. Feb 18, 2024 · This web UI, specifically designed for stable diffusion models, offers intuitive controls and options for generating text and image samples with textual inversion. Jul 6, 2024 · ComfyUI is a node-based GUI for Stable Diffusion. A model won’t be able to generate a cat’s image if there’s never a cat in the training data. 5, Stable Diffusion XL (SDXL), and Kandinsky 2. For a general introduction to the Stable Diffusion model please refer to this colab. Compared to its predecessor, the SDXL 1. Describe how the final image should look like. Starting image. Oct 25, 2022 · Training approach. The next step is to install the tools required to run stable diffusion; this step can take approximately 10 minutes. 🤗 Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. Make sure the drive you create the folder on has enough available space on it. It can be used for a wide range of applications As you’ll learn on your Stable Diffusion journey, diffusion image generation has a few weakness—especially when it comes to drawing words, symbols, and fingers. If the AI image is in PNG format, you can try to see if the prompt and other setting information were written in the PNG metadata field. Depending on the algorithm and settings, you might notice different distortions, such as gentle blurring, texture exaggeration, or color smearing. To do this, just click on the Image to Image (Img2Img) tab, place the reference image in the appropriate box, create the prompt you want the machine to follow, and click generate. This component runs for multiple steps to generate image information. Stable Diffusion image 1 using 3D rendering. Aug 4, 2023 · Image to Image essentially lets Stable Diffusion create a new image using another picture as reference, doesn’t matter whether it's a real image or one you've created. First, your text prompt gets projected into a latent vector space by the Jan 24, 2023 · Diffusion Models for Image Generation – A Comprehensive Guide. This process is similar to the diffusion process in physics, where particles spread from areas of high Sep 28, 2023 · Source: FormatPDF. Jan 7, 2023 · Stable Diffusion est une collection de modèle d'intelligence artificielle (IA) créés et partagés par Stability AI. ControlNet adds one more conditioning in addition to the text prompt. It’s where a lot of the performance gain over previous models is achieved. Set denoising strength to 0. Then run Stable Diffusion in a special python environment using Miniconda. In order to inpaint specific areas, we need to create a mask using the AUTOMATIC1111 GUI. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, cultivates autonomous freedom to produce incredible imagery, empowers billions of people to create stunning art within seconds. Jun 21, 2023 · Running the Diffusion Process. Keep reading to start creating. Note: Stable Diffusion v1 is a general text-to-image diffusion Jan 9, 2023 · Lexica is a collection of images with prompts. Img2Img is a cutting-edge technique that generates new images from an input image and a corresponding text prompt. Are you ready to revolutionize the way you work with Stable Diffusion, This is precisely the essential and comprehensive course for all of you who are passionate about Generative AI and Generating your own AI images for both your Work and Leisure. While increasing the number of steps does improve the clarity and logic of generation, it has a very sharp diminishing return curve, and rather obviously, more steps will result in longer generation times. On the Settings page, click User Interface on the left panel. Most commonly applied to reconstructing old deteriorated images, removing cracks, scratches, dust spots, or red-eyes from photographs. 75 give a good balance. We will inpaint both the right arm and the face at the same time. A text prompt. This enhances scalability, supporting models with up to 8 billion parameters and multi-modal inputs. Use the paintbrush tool to create a mask on the face. Note: Stable Diffusion v1 is a general text-to-image diffusion Oct 28, 2023 · Check out the Quick Start Guide if you are new to Stable Diffusion. In the current workflow, fixing characters and image styles often need Jul 11, 2024 · Step 3: Create a Folder for Stable Diffusion. Nov 28, 2023 · Luckily, you can use inpainting to fix it. Repeat the process until you achieve the desired outcome. Aug 28, 2023 · Navigate to the command center of Img2Img (Stable Diffusion image-to-image) – the realm where your creation takes shape. Stable Diffusion Interactive Notebook 📓 🤖. The model uses a technique called "diffusion," which generates images by gradually adding and removing noise. Early morning, with mist rising off the water, natural light, wide angle shot, shot on Canon EOS R5 with a 24mm f/11 lens” — Rob’s Mix Ultimate. You can also type in a specific seed number into this field. We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. Stable Diffusion 3 combines a diffusion transformer architecture and flow matching. Ces modèles sont utilisé pour générer des images - le plus souvent en utilisant le text-to-image : vous lui donnez un texte de description (appelé prompt en anglais) et le modèle le transforme en image. I will create it on E://. Stable Diffusion 1. The Stable Diffusion model was created by researchers and engineers from CompVis, Stability AI, Runway, and LAION. Outpainting complex scenes. a CompVis. 3 days ago · Stable Diffusion 3. The resolution has increased by 168%, from 768×768 pixels in v2 to 2048× Sep 16, 2023 · In this comprehensive guide, we’ll walk you through setting up the software, using the color sketch tool, and leveraging Img2Img to turn amateur sketches into professional artwork. A guide to Stable Diffusion . A higher value will result in more details Online. The most basic form of using Stable Diffusion models is text-to-image. Mar 13, 2024 · Stable diffusion is a deep learning model that utilizes a diffusion model to generate high-quality images based on text descriptions. For more information, you can check out Text-to-image. Values between 0. In the Stable Diffusion checkpoint dropdown menu, Select the model you originally used when generating this image . Step 3: Set outpainting parameters. Artificial Intelligence (AI) art is currently all the rage, but most AI image generators run in the cloud. extremely detailed, ornate. Stable Diffusion is an AI-powered tool that enables users to transform plain text into images. The first rule is be as specific. After applying stable diffusion, take a close look at the nudified image. Step 4: Enable the outpainting script. Negative Prompts. However, due to the granularity and method of its control, the efficiency improvement is limited for professional artistic creations such as comics and animation production whose main work is secondary painting. co, and install them. Method 1: Get prompts from images by reading PNG Info. To run Stable Diffusion locally on your PC, download Stable Diffusion from GitHub and the latest checkpoints from HuggingFace. Diffusion models, including Glide, Dalle-2, Imagen, and Stable Diffusion, have spearheaded recent advances in AI-based image generation, taking the world of “ AI Art generation ” by storm. Stable Diffusion image 2 using 3D rendering. Check out the Quick Start Guide if you are new to Stable Diffusion. Sep 25, 2022 · Stable Diffusion consists of three parts: A text encoder, which turns your prompt into a latent vector. Jun 21, 2023 · Apply the filter: Apply the stable diffusion filter to your image and observe the results. This notebook aims to be an alternative to WebUIs while offering a simple and lightweight GUI for anyone to get started Mar 20, 2023 · When a user asks Stable Diffusion to generate an output from an input image, whether that is through image-to-image (img2img) or InPaint, it initiates this process by adding noise to that input based on a seed. Stable Diffusion Web UI is a browser interface based on the Gradio library for Stable Diffusion. The stable diffusion model operates in a latent space, where it gradually destroys an image by adding noise and then reverses this process to generate a new image. It provides a user-friendly way to interact with Stable Diffusion, an open-source text-to-image generation model. This is the area you want Stable Diffusion to regenerate the image. Stable Diffusion is cool! Build Stable Diffusion “from Scratch”. It’s significantly better than previous Stable Diffusion models at incorporating text into images. In the Quicksetting List, add the following. Instead of introducing noise directly to an image, Stable Diffusion disrupts the image in the latent space with latent noise. Now you’ll see a page that looks like This only applies to image-to-image and inpainting generations. 5 Outpainting uses an approach that combines a diffusion model with an autoencoder. Sep 29, 2022 · Stable Diffusion takes two primary inputs and translates these into a fixed point in its model’s latent space: A seed integer. Nov 4, 2023 · Stable Diffusion and ControlNet have achieved excellent results in the field of image generation and synthesis. Whether you're looking for a simple inference solution or training your own diffusion models, 🤗 Diffusers is a modular toolbox that supports both. 1. Stable Diffusion is an open-source deep learning model that specializes in generating high-quality images from text descriptions. Creating an Inpaint Mask. If you’re eager to dive into the world of AI-generated art using Stable Diffusion XL, this guide will help you get started. Upload an Image. vivid. Now, input your NSFW prompts to guide the image generation process. But with the power of AI and the Stable Diffusion model, inpainting can be used to achieve more than that. If you want to run Stable Diffusion locally, you can follow these simple steps. Some commonly used blocks are Loading a Checkpoint Model, entering a prompt, specifying a sampler, etc. Jul 9, 2023 · 1. Aug 7, 2023 · This dataset includes images of all different sizes and shapes, so Stable Diffusion knows how to generate new pixels that match the style and content of the original image. Center an image. The dice button to the right of the Seed field will reset it to -1. It involves inputting data, allowing the AI to process through Gaussian noise, and receiving an artistic output. The model has been trained on billions of images and can produce images that are highly coherent to the ones obtainable from DALL-E 2 and MidJourney. Using Stable Diffusion is fundamentally straightforward. Stable Diffusion - InvokeAI: Supports the most features, but struggles with 4 GB or less VRAM, requires an Nvidia GPU; Stable Diffusion - OptimizedSD: Lacks many features, but runs on 4 GB or even less VRAM, requires an Nvidia GPU Apr 16, 2024 · A black and white image is used as a mask for inpainting over the provided image. The amount of noise it adds is controlled by Denoising Strength, which can be a minimum of 0 and a maximum of 1. k. How strongly the original image should be altered (from subtle to drastic changes) 80 %. Inpainting & Outpainting. To produce an image, Stable Diffusion first generates a completely random image in the latent space. The StableDiffusionPipeline is capable of generating photorealistic images given any text input. You can use this software on Windows, Mac, or Google Colab. By default, it’s set to 32 pixels. May 12, 2023 · The image and prompts will be populated automatically. What kind of images a model generates depends on the training images. Aug 3, 2023 · Undoubtedly, the emergence of Stable Diffusion XL has marked a milestone in the history of natural language processing and image generation, taking us a step closer to something that already scares… Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. pruned. It’s trained on 512x512 images from a subset of the LAION-5B dataset. In this case, images Aug 30, 2023 · Diffusion Explainer is a perfect tool for you to understand Stable Diffusion, a text-to-image model that transforms a text prompt into a high-resolution image. The subject’s images are fitted alongside images from the subject’s class, which are first generated using the same Stable Diffusion model. ComfyUI breaks down a workflow into rearrangeable elements so you can easily make your own. Copy the prompt, paste it to the Stable Diffusion and press Generate to see generated images. Remember, you have the freedom to experiment with other models as well. Step 8: Generate NSFW Images. The model then creates a random image and gradually refines it by making small changes at each step. Upload the original image to be modified. Sep 22, 2022 · Wondering how to generate NSFW images in Stable Diffusion? We will show you, so you don't need to worry about filters or censorship. Failure example of Stable Diffusion outpainting. Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION. We'll talk about txt2img, img2img, Step 5: Setup the Web-UI. Dec 5, 2023 · Stable Diffusion is a text-to-image model powered by AI that can create images from text and in this guide, I'll cover all the basics. Open your command prompt and navigate to the stable-diffusion-webui folder using the following command: cd path / to / stable - diffusion - webui. 3. Popular models. Generating high-quality images from text descriptions is a challenging task. By leveraging prompt template files, users can quickly configure the web UI to generate text that aligns with specific concepts. Generate NSFW Now. Prompt: A beautiful ((Ukrainian Girl)) with very long straight hair, full lips, a gentle look, and very light white skin. Oct 10, 2023 · Stable Diffusion XL (SDXL) is a groundbreaking text-to-image generation model developed by Stability AI. There are two primary types of masks used in this process: Mask and Invert Mask. Le champs de texte Prompt (press Ctrl+Enter or All+Enter to generate) est le seul prompt requis pour pouvoir créer une image. Nov 12, 2022 · Two advices: (1) Be as detailed and specific, and (2) use powerful keywords. 3. powered by Stable Diffusion AI ( CreativeML Open RAIL-M) Prompt. You can construct an image generation workflow by chaining different blocks (called nodes) together. Higher numbers change more of the image, lower numbers keep the original image intact. She wears a medieval dress. Mask: This is used to specify the areas in an Keep in mind that, even if you don’t save the image with this method, it is still saved in an output directory specified by Stable Diffusion. zr cf ce wn pz qp kx ek de lf